Free trial included. Once the download is complete, navigate to the file on your computer and double-click to begin the installation process. I want to start by saying thank you to everyone who made Stable Diffusion UI possible. Evaluation. First experiments with SXDL, part III: Model portrait shots in automatic 1111. Only Nvidia cards are officially supported. To run Stable Diffusion via DreamStudio: Navigate to the DreamStudio website. With a ControlNet model, you can provide an additional control image to condition and control Stable Diffusion generation. Resumed for another 140k steps on 768x768 images. Remove objects, people, text and defects from your pictures automatically. This began as a personal collection of styles and notes. Ultrafast 10 Steps Generation!! (one second. いま一部で話題の Stable Diffusion 。. License: SDXL 0. ckpt file contains the entire model and is typically several GBs in size. 5 and 2. Overview. Unlike models like DALL. Step 5: Launch Stable Diffusion. This neg embed isn't suited for grim&dark images. In the folder navigate to models » stable-diffusion and paste your file there. weight, lora_down. 8 or later on your computer to run Stable Diffusion. This video is 2160x4096 and 33 seconds long. 12. Stable Diffusion XL (SDXL), is the latest AI image generation model that can generate realistic faces, legible text within the images, and better image composition, all while using shorter and simpler prompts. Model type: Diffusion-based text-to-image generative model. Cmdr2's Stable Diffusion UI v2. 0 base model & LORA: – Head over to the model. Does anyone knows if is a issue on my end or. I was looking at that figuring out all the argparse commands. I appreciate all the good feedback from the community. Stable Diffusion + ControlNet. The GPUs required to run these AI models can easily. April 11, 2023. Stable Diffusion x2 latent upscaler model card. 1. Credit: ai_coo#2852 (street art) Stable Diffusion embodies the best features of the AI art world: it’s arguably the best existing AI art model and open source. 𝑡→ 𝑡−1 •Score model 𝜃: ×0,1→ •A time dependent vector field over space. Another experimental VAE made using the Blessed script. With 256x256 it was on average 14s/iteration, so much more reasonable, but still sluggish af. 使用stable diffusion制作多人图。. Click on Command Prompt. Training diffusion model = Learning to denoise •If we can learn a score model 𝜃 , ≈∇log ( , ) •Then we can denoise samples, by running the reverse diffusion equation. I wanted to document the steps required to run your own model and share some tips to ensure that you are starting on the right foot. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. Notice there are cases where the output is barely recognizable as a rabbit. 如果想要修改. Stable Diffusion is a deep learning based, text-to-image model. Compared to. patrickvonplaten HF staff. • 4 mo. Use your browser to go to the Stable Diffusion Online site and click the button that says Get started for free. Sounds Like a Metal Band: Fun with DALL-E and Stable Diffusion. Look at the file links at. Keyframes created and link to method in the first comment. A dmg file should be downloaded. paths import script_path line after from. Reload to refresh your session. 0 is released. 5. Bryan Bischof Sep 8 GenAI, Stable Diffusion, DALL-E, Computer. com models though are heavily scewered in specific directions, if it comes to something that isnt anime, female pictures, RPG, and a few other popular themes then it still performs fairly poorly. r/StableDiffusion. Stability AI Ltd. py", line 185, in load_lora assert False, f'Bad Lora layer name: {key_diffusers} - must end in lora_up. The difference is subtle, but noticeable. fix to scale it to whatever size I want. Paper: "Beyond Surface Statistics: Scene Representations in a Latent Diffusion Model". Thanks to a generous compute donation from Stability AI and support from LAION, we were able to train a Latent Diffusion Model on 512x512 images from a subset of the LAION-5B database. Artist Inspired Styles. It helps blend styles together! 1 / 7. Our Language researchers innovate rapidly and release open models that rank amongst the best in the. It can be used in combination with Stable Diffusion. Additional training is achieved by training a base model with an additional dataset you are. Stable Diffusion XL lets you create better, bigger pictures, with faces that look more real. Height. This checkpoint corresponds to the ControlNet conditioned on HED Boundary. Generate the image. Includes the ability to add favorites. 手順2：「gui. You can disable hardware acceleration in the Chrome settings to stop it from using any VRAM, will help a lot for stable diffusion. 0. The GPUs required to run these AI models can easily. Step. 10. Stable Diffusion XL 1. Appendix A: Stable Diffusion Prompt Guide. 1, which both failed to replace their predecessor. SToday, Stability AI announces SDXL 0. SDXL 0. Those will probably be need to be fed to the 'G' Clip of the text encoder. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet. Log in. Experience cutting edge open access language models. → Stable Diffusion v1モデル_H2. Stable Diffusion XL delivers more photorealistic results and a bit of text In general, SDXL seems to deliver more accurate and higher quality results, especially in the area of photorealism. from_pretrained( "stabilityai/stable-diffusion. With ComfyUI it generates images with no issues, but it's about 5x slower overall than SD1. 5 and 2. . It is common to see extra or missing limbs. We are building the foundation to activate humanity's potential. License: SDXL 0. Try Stable Audio Stable LM. Select “stable-diffusion-v1-4. If a seed is provided, the resulting. b) for sanity check, i would try the LoRA model on a painting/illustration focused stable diffusion model (anime checkpoints works) and see if the face is recognizable, if it is, it is an indication to me that the LoRA is trained "enough" and the concept should be transferable for most of my use. No VAE compared to NAI Blessed. In this video, I will show you how to install **Stable Diffusion XL 1. While this model hit some of the key goals I was reaching for, it will continue to be trained to fix the weaknesses. It'll always crank up the exposure and saturation or neglect prompts for dark exposure. Stability AI. Download the SDXL 1. History: 18 commits. Human anatomy, which even Midjourney struggled with for a long time, is also handled much better by SDXL, although the finger problem seems to have. Head to Clipdrop, and select Stable Diffusion XL (or just click here ). Stable Diffusion XL (SDXL) was proposed in SDXL: Improving Latent Diffusion Models for High-Resolution Image Synthesis by Dustin Podell, Zion English, Kyle Lacey, Andreas Blattmann, Tim Dockhorn, Jonas Müller, Joe Penna, and Robin Rombach. "art in the style of Amanda Sage" 40 steps. Follow the link below to learn more and get installation instructions. Note that it will return a black image and a NSFW boolean. Stable Diffusion is a Latent Diffusion model developed by researchers from the Machine Vision and Learning group at LMU Munich, a. Stable Diffusion. that slows down stable diffusion. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. Try on Clipdrop. • 19 days ago. 9 and Stable Diffusion 1. Now you can set any count of images and Colab will generate as many as you set On Windows - WIP Prerequisites . Loading weights [5c5661de] from D:AIstable-diffusion-webuimodelsStable-diffusionx4-upscaler-ema. I like how you have put a different prompt into your upscaler and ControlNet than the main prompt: I think this could help to stop getting random heads from appearing in tiled upscales. ago. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. 9. A text-to-image generative AI model that creates beautiful images. SD-XL. k. Especially on faces. Some types of picture include digital illustration, oil painting (usually good results), matte painting, 3d render, medieval map. 0 Model. 0, an open model representing the next evolutionary step in text-to. Forward diffusion gradually adds noise to images. Steps. First, the stable diffusion model takes both a latent seed and a text prompt as input. You'll see this on the txt2img tab:I made a long guide called [Insights for Intermediates] - How to craft the images you want with A1111, on Civitai. 0 and try it out for yourself at the links below : SDXL 1. Quick Tip for Beginners: You can change the default settings of Stable Diffusion WebUI (AUTOMATIC1111) in the ui-config. md. today introduced Stable Audio, a software platform that uses a latent diffusion model to generate audio based on users' text prompts. 本教程需要一些AI绘画基础，并不是面对0基础人员，如果你没有学习过stable diffusion的基本操作或者对Controlnet插件毫无了解，可以先看看秋葉aaaki等up的教程，做到会存放大模型，会安装插件并且有基本的视频剪辑能力。-----一、准备工作Launching Web UI with arguments: --xformers Loading weights [dcd690123c] from C: U sers d alto s table-diffusion-webui m odels S table-diffusion v 2-1_768-ema-pruned. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input. AUTOMATIC1111 / stable-diffusion-webui. py ", line 294, in lora_apply_weights. This checkpoint corresponds to the ControlNet conditioned on M-LSD straight line detection. Stable Diffusion is one of the most famous examples that got wide adoption in the community and. true. Comparison. First, describe what you want, and Clipdrop Stable Diffusion XL will generate four pictures for you. . The structure of the prompt. 225,000 steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10 % dropping of the text-conditioning to improve classifier-free guidance sampling. . I created a trailer for a Lakemonster movie with MidJourney, Stable Diffusion and other AI tools. Enter a prompt, and click generate. “The audio quality is astonishing. 5. 手順3：学習を行う. Additionally, their formulation allows for a guiding mechanism to control the image generation process without retraining. diffusion_pytorch_model. Figure 3: Latent Diffusion Model (Base Diagram:[3], Concept-Map Overlay: Author) A very recent proposed method which leverages upon the perceptual power of GANs, the detail preservation ability of the Diffusion Models, and the Semantic ability of Transformers by merging all three together. My A1111 takes FOREVER to start or to switch between checkpoints because it's stuck on "Loading weights [31e35c80fc] from a1111stable-diffusion-webuimodelsStable-diffusionsd_xl_base_1. Diffusion models are a. invokeai is always a good option. yaml LatentUpscaleDiffusion: Running in v-prediction mode DiffusionWrapper has 473. . Learn. Latent Diffusion models are game changers when it comes to solving text-to-image generation problems. Step 2: Double-click to run the downloaded dmg file in Finder. Check out my latest video showing Stable Diffusion SXDL for hi-res AI… Liked by Oliver Hamilton. Below are some of the key features: – User-friendly interface, easy to use right in the browser – Supports various image generation options like size, amount, mode,. Skip to main contentModel type: Diffusion-based text-to-image generative model. SDXL - The Best Open Source Image Model. In this blog post, we will: Explain the. However, since these models. 1 embeddings, hypernetworks and Loras. Sort by: Open comment sort options. proj_in in the given object!. 0からは花札アイコンは消えてデフォルトでタブ表示になりました。Stable diffusion 配合 ControlNet 骨架分析，输出的图片确实让人大吃一惊！. This model runs on Nvidia A40 (Large) GPU hardware. I said earlier that a prompt needs to be detailed and specific. civitai. 0 (Stable Diffusion XL) has been released earlier this week which means you can run the model on your own computer and generate images using your own GPU. Stable Diffusion Desktop Client. I am pleased to see the SDXL Beta model has. . Your image will be generated within 5 seconds. 前提：Stable. Download the Latest Checkpoint for Stable Diffusion from Huggin Face. Click to open Colab link . It serves as a quick reference as to what the artist's style yields. By default, Colab notebooks rely on the original Stable Diffusion which comes with NSFW filters. ️ Check out Lambda here and sign up for their GPU Cloud: it here online: to run it:. DiffusionBee is one of the easiest ways to run Stable Diffusion on Mac. Stable Diffusion creates an image by starting with a canvas full of noise and denoise it gradually to reach the final output. The Stable-Diffusion-v1-5 checkpoint was initialized with the weights of the Stable-Diffusion-v1-2 checkpoint and subsequently fine-tuned on 595k steps at resolution 512x512 on "laion-aesthetics v2 5+" and 10%. Model Description: This is a model that can be used to generate and modify images based on text prompts. Not a LORA, but you can download ComfyUI nodes for sharpness, blur, contrast, saturation, sharpness, etc. In general, the best stable diffusion prompts will have this form: “A [type of picture] of a [main subject], [style cues]* ”. Those will probably be need to be fed to the 'G' Clip of the text encoder. We’re on a journey to advance and democratize artificial intelligence through. For more information, you can check out. 9, a follow-up to Stable Diffusion XL. ckpt - format is commonly used to store and save models. Stable Diffusion 1. 1. 9 Tutorial (better than Midjourney AI)Stability AI recently released SDXL 0. Model type: Diffusion-based text-to-image generation modelStable Diffusion XL. Stable Diffusion is a deep learning generative AI model. Paper: "Beyond Surface Statistics: Scene Representations in a Latent Diffusion Model". 0, the flagship image model developed by Stability AI, stands as the pinnacle of open models for image generation. Specializing in ultra-high-resolution outputs, it's the ideal tool for producing large-scale artworks and. With 3. # How to turn a painting into a landscape via SXDL Controlnet ComfyUI: 1. SDXL is supposedly better at generating text, too, a task that’s historically. Its the guide that I wished existed when I was no longer a beginner Stable Diffusion user. 【Stable Diffusion】 超强AI绘画，FeiArt教你在线免费玩！想深入探讨，可以加入FeiArt创建的AI绘画交流扣扣群：926267297我们在群里目前搭建了免费的国产Ai绘画机器人，大家可以直接试用。后续可能也会搭建SD版本的绘画机器人群。免费在线体验Stable diffusion链接：无需注册和充钱版，但要排队：. The Stable Diffusion 1. Download the zip file and use it as your own personal cheat-sheet - completely offline. Tried with a base model 8gb m1 mac. Can someone for the love of whoever is most dearest to you post a simple instruction where to put the SDXL files and how to run the thing?. Then you can pass a prompt and the image to the pipeline to generate a new image:Summary. Easy Diffusion is a simple way to download Stable Diffusion and use it on your computer. In this post, you will learn the mechanics of generating photo-style portrait images. seed: 1. By simply replacing all instances linking to the original script with the script that has no safety filters, you can easily achieve generate NSFW images. SDXL consists of an ensemble of experts pipeline for latent diffusion: In a first step, the base model is used to generate (noisy) latents, which are then further processed with a refinement model (available here: specialized for the final denoising steps. 9, which adds image-to-image generation and other capabilities. Temporalnet is a controlNET model that essentially allows for frame by frame optical flow, thereby making video generations significantly more temporally coherent. Understandable, it was just my assumption from discussions that the main positive prompt was for common language such as "beautiful woman walking down the street in the rain, a large city in the background, photographed by PhotographerName" and the POS_L and POS_R would be for detailing such as "hyperdetailed, sharp focus, 8K, UHD" that sort of thing. 9 Tutorial (better than Midjourney AI)Stability AI recently released SDXL 0. Copy the file, and navigate to Stable Diffusion folder you created earlier. Download all models and put into stable-diffusion-webuimodelsStable-diffusion folder; Test with run. afaik its only available for inside commercial teseters presently. When conducting densely conditioned tasks with the model, such as super-resolution, inpainting, and semantic synthesis, the stable diffusion model is able to generate megapixel images (around 10242 pixels in size). Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. Stable Diffusion is the latest deep learning model to generate brilliant, eye-catching art based on simple input text. fp16. 0 model. On the one hand it avoids the flood of nsfw models from SD1. Waiting at least 40s per generation (comfy; the best performance I've had) is tedious and I don't have much free. KOHYA. Developed by: Stability AI. Checkpoints, Loras, hypernetworks, text inversions, and prompt words. 8 GB LoRA Training - Fix CUDA Version For DreamBooth and Textual Inversion Training By Automatic1111. 9 and Stable Diffusion 1. 0. 5 I used Dreamshaper 6 since it's one of the most popular and versatile models. The platform can generate up to 95-second cli，相关视频：sadtalker安装中的疑难杂症帮你搞定，SadTalker最新版本安装过程详解，Stable Diffusion 一键安装包，秋叶安装包，AI安装包，一键部署，stable diffusion 秋叶4. Try to reduce those to the best 400 if you want to capture the style. The only caveat here is that you need a Colab Pro account since the free version of Colab offers not enough VRAM to. 147. Click to see where Colab generated images. ago. 5d4cfe8 about 1 month ago. We're going to create a folder named "stable-diffusion" using the command line. You can try it out online at beta. While these are not the only solutions, these are accessible and feature rich, able to support interests from the AI art-curious to AI code warriors. 0 with the current state of SD1. The late-stage decision to push back the launch "for a week or so," disclosed by Stability AI’s Joe. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. At the time of release (October 2022), it was a massive improvement over other anime models. 5 base. :( Almost crashed my PC! Stable LM. Model Details Developed by: Lvmin Zhang, Maneesh Agrawala. 5. You’ll also want to make sure you have 16 GB of PC RAM in the PC system to avoid any instability. By decomposing the image formation process into a sequential application of denoising autoencoders, diffusion models (DMs) achieve state-of-the-art synthesis results on image data and beyond. Join. Copy and paste the code block below into the Miniconda3 window, then press Enter. Learn more about A1111. 0 and was released in lllyasviel/ControlNet-v1-1 by Lvmin Zhang. Use "Cute grey cats" as your prompt instead. With Stable Diffusion XL, you can create descriptive images with shorter prompts and generate words within images. torch. 4版本+WEBUI1. Unsupervised Semantic Correspondences with Stable Diffusion to appear at NeurIPS 2023. Fine-tuned Model Checkpoints (Dreambooth Models) Download the custom model in Checkpoint format (. First create a new conda environmentLearn more about Stable Diffusion SDXL 1. Chrome uses a significant amount of VRAM. You switched accounts on another tab or window. #SDXL is currently in beta and in this video I will show you how to use it on Google Colab for free. It. 1 is the successor model of Controlnet v1. 1. It can generate novel images from text descriptions and produces. There is still room for further growth compared to the improved quality in generation of hands. safetensors Creating model from config: C: U sers d alto s table-diffusion-webui epositories s table-diffusion-stability-ai c onfigs s table-diffusion v 2-inference. 0: cfg_scale – How strictly the diffusion process adheres to the prompt text. 1. This base model is available for download from the Stable Diffusion Art website. Run the command conda env create -f environment. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. Just like its. You signed in with another tab or window. The world of AI image generation has just taken another significant leap forward. 0, an open model representing the next evolutionary step in text-to-image generation models. This version of Stable Diffusion creates a server on your local PC that is accessible via its own IP address, but only if you connect through the correct port: 7860. TypeScript. 2, along with code to get started with deploying to Apple Silicon devices. Some types of picture include digital illustration, oil painting (usually good results), matte painting, 3d render, medieval map. attentions. 9 base model gives me much(!) better results with the. Create a folder in the root of any drive (e. In general, the best stable diffusion prompts will have this form: “A [type of picture] of a [main subject], [style cues]* ”. They both start with a base model like Stable Diffusion v1. It's trained on 512x512 images from a subset of the LAION-5B database. Type cmd. I personally prefer 0. com不然我骚扰你. 0 + Automatic1111 Stable Diffusion webui. In a groundbreaking announcement, Stability AI has unveiled SDXL 0. Update README. Replicate was ready from day one with a hosted version of SDXL that you can run from the web or using our cloud API. Posted by 13 hours ago. Dreambooth is considered more powerful because it fine-tunes the weight of the whole model. Note: Earlier guides will say your VAE filename has to have the same as your model. compile support. 0 (SDXL 1. stable. Below are three emerging solutions for doing Stable Diffusion Generative AI art using Intel Arc GPUs on a Windows laptop or PC. Model Access Each checkpoint can be used both with Hugging Face's 🧨 Diffusers library or the original Stable Diffusion GitHub repository. Welcome to Stable Diffusion; the home of Stable Models and the Official Stability. 人物面部、手部，及背景的任意替换，手部修复的替代办法，Segment Anything +ControlNet 的灵活应用,念咒结束，【入门02】喂饭级stable diffusion安装教程，有手就会. 0 should be placed in a directory. The formula is this (epochs are useful so you can test different loras outputs per epoch if you set it like that): [ [images] x [repeats]] x [epochs] / [batch] = [total steps] Nezarah. the SXDL doesn't bring anything new to the table, maybe 0. Once you are in, input your text into the textbox at the bottom, next to the Dream button. Enter a prompt and a URL to generate. This step downloads the Stable Diffusion software (AUTOMATIC1111). from diffusers import StableDiffusionXLPipeline, StableDiffusionXLImg2ImgPipeline import torch pipeline = StableDiffusionXLPipeline. This video is 2160x4096 and 33 seconds long. You can modify it, build things with it and use it commercially. DreamStudioのアカウント作成. 6 API acts as a replacement for Stable Diffusion 1. Use it with 🧨 diffusers. It is a diffusion model that operates in the same latent space as the Stable Diffusion model. Unconditional image generation Text-to-image Stable Diffusion XL Kandinsky 2. 368. LAION-5B is the largest, freely accessible multi-modal dataset that currently exists. SDXL 0. Alternatively, you can access Stable Diffusion non-locally via Google Colab. The secret sauce of Stable Diffusion is that it "de-noises" this image to look like things we know about. For SD1. It is the best multi-purpose. PARASOL GIRL. The original Stable Diffusion model was created in a collaboration with CompVis and RunwayML and builds upon the work: High-Resolution Image Synthesis with Latent Diffusion Models. Now Stable Diffusion returns all grey cats. 这可能是唯一能将stable diffusion讲清楚的教程了，不愧是腾讯大佬！ 1天全面了解stable diffusion最全使用攻略！ ，AI绘画基础-01Stable Diffusion零基础入门，2023年11月版最新版ChatGPT和GPT 4. kohya SS gui optimal parameters - Kohya DyLoRA , Kohya LoCon , LyCORIS/LoCon , LyCORIS/LoHa , Standard Question | Helpfast-stable-diffusion Notebooks, A1111 + ComfyUI + DreamBooth. I have been using Stable Diffusion UI for a bit now thanks to its easy Install and ease of use, since I had no idea what to do or how stuff works. com はじめに今回の学習は「DreamBooth fine-tuning of the SDXL UNet via LoRA」として紹介されています。いわゆる通常のLoRAとは異なるようです。16GBで動かせるということはGoogle Colabで動かせるという事だと思います。自分は宝の持ち腐れのRTX 4090をここぞとばかりに使いました。 touch-sp. The most important shift that Stable Diffusion 2 makes is replacing the text encoder. 35. I would hate to start from zero again. This ability emerged during the training phase of the AI, and was not programmed by people. Stable Diffusion is a text-to-image latent diffusion model created by the researchers and engineers from CompVis, Stability AI and LAION. 12 Keyframes, all created in Stable Diffusion with temporal consistency. 4发. You will see the exact keyword applied to two classes of images: (1) a portrait and (2) a scene. 368. Intel Arc A750 and A770 review: Trouncing NVIDIA and AMD on mid-range GPU value | Engadget engadget. Stable Diffusion XL (SDXL) is the new open-source image generation model created by Stability AI that represents a major advancement in AI text-to-image technology. g. safetensors" I dread every time I have to restart the UI. 5. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Anyways those are my initial impressions!. down_blocks. Useful support words: excessive energy, scifi Original SD1. It is primarily used to generate detailed images conditioned on text descriptions. 0 base specifically. 4发布！ How to Train a Stable Diffusion Model Stable diffusion technology has emerged as a game-changer in the field of artificial intelligence, revolutionizing the way models are… 8 min read · Jul 18 Start stable diffusion; Choose Model; Input prompts, set size, choose steps (doesn't matter how many, but maybe with fewer steps the problem is worse), cfg scale doesn't matter too much (within limits) Run the generation; look at the output with step by step preview on. Stable Diffusion XL. 本日、 Stability AIは、フォトリアリズムに優れたエンタープライズ向け最新画像生成モデル「Stabile Diffusion XL（SDXL）」をリリースしたことを発表しました。 SDXLは、エンタープライズ向けにStability AIのAPIを通じて提供されるStable Diffusion のモデル群に新たに追加されたものです。This is an answer that someone corrects. 9. 2022/08/27. 9, which. You will usually use inpainting to correct them. Think of them as documents that allow you to write and execute code all. Model 1. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders ( OpenCLIP-ViT/G and CLIP-ViT/L ). 5; DreamShaper; Kandinsky-2; DeepFloyd IF. We’re on a journey to advance and democratize artificial intelligence through open source and open science. Hot New Top. 6 API!This API is designed to be a higher quality, more cost-effective alternative to stable-diffusion-v1-5 and is ideal for users who are looking to replace it in their workflows. High resolution inpainting - Source. A group of open source hackers forked Stable Diffusion on GitHub and optimized the model to run on Apple's M1 chip, enabling images to be generated in ~ 15 seconds (512x512 pixels, 50 diffusion steps).